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package/docs/80_upstream-docs-and-skills/hermes-dev-skill/src/i18n/zh-Hans/user-guide/skills/optional/mlops/mlops-stable-diffusion.md

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----
-title: "Stable Diffusion 图像生成"
-sidebar_label: "Stable Diffusion 图像生成"
-description: "通过 HuggingFace Diffusers 使用 Stable Diffusion 模型实现最先进的文本到图像生成"
----
-
-{/* This page is auto-generated from the skill's SKILL.md by website/scripts/generate-skill-docs.py. Edit the source SKILL.md, not this page. */}
-
-# Stable Diffusion 图像生成
-
-通过 HuggingFace Diffusers 使用 Stable Diffusion 模型实现最先进的文本到图像生成。适用于从文本 prompt（提示词）生成图像、执行图像到图像转换、图像修复（inpainting），或构建自定义扩散 pipeline。
-
-## Skill 元数据
-
-| | |
-|---|---|
-| 来源 | 可选 — 通过 `hermes skills install official/mlops/stable-diffusion` 安装 |
-| 路径 | `optional-skills/mlops/stable-diffusion` |
-| 版本 | `1.0.0` |
-| 作者 | Orchestra Research |
-| 许可证 | MIT |
-| 依赖项 | `diffusers>=0.30.0`, `transformers>=4.41.0`, `accelerate>=0.31.0`, `torch>=2.0.0` |
-| 平台 | linux, macos, windows |
-| 标签 | `Image Generation`, `Stable Diffusion`, `Diffusers`, `Text-to-Image`, `Multimodal`, `Computer Vision` |
-
-## 参考：完整 SKILL.md
-
-:::info
-以下是 Hermes 在触发此 skill 时加载的完整 skill 定义。这是 agent 在 skill 激活时所看到的指令内容。
-:::
-
-# Stable Diffusion 图像生成
-
-使用 HuggingFace Diffusers 库通过 Stable Diffusion 生成图像的综合指南。
-
-## 何时使用 Stable Diffusion
-
-**在以下情况下使用 Stable Diffusion：**
-- 从文本描述生成图像
-- 执行图像到图像转换（风格迁移、增强）
-- Inpainting（填充遮罩区域）
-- Outpainting（将图像扩展至边界之外）
-- 创建现有图像的变体
-- 构建自定义图像生成工作流
-
-**核心功能：**
-- **文本到图像**：从自然语言 prompt 生成图像
-- **图像到图像**：在文本引导下转换现有图像
-- **Inpainting**：用上下文感知内容填充遮罩区域
-- **ControlNet**：添加空间条件控制（边缘、姿态、深度）
-- **LoRA 支持**：高效微调与风格适配
-- **多模型支持**：支持 SD 1.5、SDXL、SD 3.0、Flux
-
-**改用以下替代方案：**
-- **DALL-E 3**：无需 GPU 的 API 生成
-- **Midjourney**：艺术化、风格化输出
-- **Imagen**：Google Cloud 集成
-- **Leonardo.ai**：基于 Web 的创意工作流
-
-## 快速开始
-
-### 安装
-
-```bash
-pip install diffusers transformers accelerate torch
-pip install xformers  # Optional: memory-efficient attention
-```
-
-### 基础文本到图像
-
-```python
-from diffusers import DiffusionPipeline
-import torch
-
-# Load pipeline (auto-detects model type)
-pipe = DiffusionPipeline.from_pretrained(
-    "stable-diffusion-v1-5/stable-diffusion-v1-5",
-    torch_dtype=torch.float16
-)
-pipe.to("cuda")
-
-# Generate image
-image = pipe(
-    "A serene mountain landscape at sunset, highly detailed",
-    num_inference_steps=50,
-    guidance_scale=7.5
-).images[0]
-
-image.save("output.png")
-```
-
-### 使用 SDXL（更高质量）
-
-```python
-from diffusers import AutoPipelineForText2Image
-import torch
-
-pipe = AutoPipelineForText2Image.from_pretrained(
-    "stabilityai/stable-diffusion-xl-base-1.0",
-    torch_dtype=torch.float16,
-    variant="fp16"
-)
-pipe.to("cuda")
-
-# Enable memory optimization
-pipe.enable_model_cpu_offload()
-
-image = pipe(
-    prompt="A futuristic city with flying cars, cinematic lighting",
-    height=1024,
-    width=1024,
-    num_inference_steps=30
-).images[0]
-```
-
-## 架构概览
-
-### 三支柱设计
-
-Diffusers 围绕三个核心组件构建：
-
-<!-- ascii-guard-ignore -->
-```
-Pipeline (orchestration)
-├── Model (neural networks)
-│   ├── UNet / Transformer (noise prediction)
-│   ├── VAE (latent encoding/decoding)
-│   └── Text Encoder (CLIP/T5)
-└── Scheduler (denoising algorithm)
-```
-<!-- ascii-guard-ignore-end -->
-
-### Pipeline 推理流程
-
-```
-Text Prompt → Text Encoder → Text Embeddings
-                                    ↓
-Random Noise → [Denoising Loop] ← Scheduler
-                      ↓
-               Predicted Noise
-                      ↓
-              VAE Decoder → Final Image
-```
-
-## 核心概念
-
-### Pipeline
-
-Pipeline 编排完整工作流：
-
-| Pipeline | 用途 |
-|----------|---------|
-| `StableDiffusionPipeline` | 文本到图像（SD 1.x/2.x） |
-| `StableDiffusionXLPipeline` | 文本到图像（SDXL） |
-| `StableDiffusion3Pipeline` | 文本到图像（SD 3.0） |
-| `FluxPipeline` | 文本到图像（Flux 模型） |
-| `StableDiffusionImg2ImgPipeline` | 图像到图像 |
-| `StableDiffusionInpaintPipeline` | Inpainting |
-
-### Scheduler
-
-Scheduler 控制去噪过程：
-
-| Scheduler | 步数 | 质量 | 适用场景 |
-|-----------|-------|---------|----------|
-| `EulerDiscreteScheduler` | 20-50 | 良好 | 默认选择 |
-| `EulerAncestralDiscreteScheduler` | 20-50 | 良好 | 更多变化 |
-| `DPMSolverMultistepScheduler` | 15-25 | 优秀 | 快速、高质量 |
-| `DDIMScheduler` | 50-100 | 良好 | 确定性生成 |
-| `LCMScheduler` | 4-8 | 良好 | 极速生成 |
-| `UniPCMultistepScheduler` | 15-25 | 优秀 | 快速收敛 |
-
-### 切换 Scheduler
-
-```python
-from diffusers import DPMSolverMultistepScheduler
-
-# Swap for faster generation
-pipe.scheduler = DPMSolverMultistepScheduler.from_config(
-    pipe.scheduler.config
-)
-
-# Now generate with fewer steps
-image = pipe(prompt, num_inference_steps=20).images[0]
-```
-
-## 生成参数
-
-### 关键参数
-
-| 参数 | 默认值 | 说明 |
-|-----------|---------|-------------|
-| `prompt` | 必填 | 目标图像的文本描述 |
-| `negative_prompt` | None | 图像中需要避免的内容 |
-| `num_inference_steps` | 50 | 去噪步数（越多质量越好） |
-| `guidance_scale` | 7.5 | Prompt 遵循程度（通常为 7-12） |
-| `height`, `width` | 512/1024 | 输出尺寸（8 的倍数） |
-| `generator` | None | 用于可复现性的 Torch generator |
-| `num_images_per_prompt` | 1 | 批量大小 |
-
-### 可复现生成
-
-```python
-import torch
-
-generator = torch.Generator(device="cuda").manual_seed(42)
-
-image = pipe(
-    prompt="A cat wearing a top hat",
-    generator=generator,
-    num_inference_steps=50
-).images[0]
-```
-
-### Negative prompt
-
-```python
-image = pipe(
-    prompt="Professional photo of a dog in a garden",
-    negative_prompt="blurry, low quality, distorted, ugly, bad anatomy",
-    guidance_scale=7.5
-).images[0]
-```
-
-## 图像到图像
-
-在文本引导下转换现有图像：
-
-```python
-from diffusers import AutoPipelineForImage2Image
-from PIL import Image
-
-pipe = AutoPipelineForImage2Image.from_pretrained(
-    "stable-diffusion-v1-5/stable-diffusion-v1-5",
-    torch_dtype=torch.float16
-).to("cuda")
-
-init_image = Image.open("input.jpg").resize((512, 512))
-
-image = pipe(
-    prompt="A watercolor painting of the scene",
-    image=init_image,
-    strength=0.75,  # How much to transform (0-1)
-    num_inference_steps=50
-).images[0]
-```
-
-## Inpainting
-
-填充遮罩区域：
-
-```python
-from diffusers import AutoPipelineForInpainting
-from PIL import Image
-
-pipe = AutoPipelineForInpainting.from_pretrained(
-    "runwayml/stable-diffusion-inpainting",
-    torch_dtype=torch.float16
-).to("cuda")
-
-image = Image.open("photo.jpg")
-mask = Image.open("mask.png")  # White = inpaint region
-
-result = pipe(
-    prompt="A red car parked on the street",
-    image=image,
-    mask_image=mask,
-    num_inference_steps=50
-).images[0]
-```
-
-## ControlNet
-
-添加空间条件控制以实现精确控制：
-
-```python
-from diffusers import StableDiffusionControlNetPipeline, ControlNetModel
-import torch
-
-# Load ControlNet for edge conditioning
-controlnet = ControlNetModel.from_pretrained(
-    "lllyasviel/control_v11p_sd15_canny",
-    torch_dtype=torch.float16
-)
-
-pipe = StableDiffusionControlNetPipeline.from_pretrained(
-    "stable-diffusion-v1-5/stable-diffusion-v1-5",
-    controlnet=controlnet,
-    torch_dtype=torch.float16
-).to("cuda")
-
-# Use Canny edge image as control
-control_image = get_canny_image(input_image)
-
-image = pipe(
-    prompt="A beautiful house in the style of Van Gogh",
-    image=control_image,
-    num_inference_steps=30
-).images[0]
-```
-
-### 可用的 ControlNet
-
-| ControlNet | 输入类型 | 适用场景 |
-|------------|------------|----------|
-| `canny` | 边缘图 | 保留结构 |
-| `openpose` | 姿态骨架 | 人体姿态 |
-| `depth` | 深度图 | 3D 感知生成 |
-| `normal` | 法线图 | 表面细节 |
-| `mlsd` | 线段 | 建筑线条 |
-| `scribble` | 粗略草图 | 草图到图像 |
-
-## LoRA 适配器
-
-加载微调风格适配器：
-
-```python
-from diffusers import DiffusionPipeline
-
-pipe = DiffusionPipeline.from_pretrained(
-    "stable-diffusion-v1-5/stable-diffusion-v1-5",
-    torch_dtype=torch.float16
-).to("cuda")
-
-# Load LoRA weights
-pipe.load_lora_weights("path/to/lora", weight_name="style.safetensors")
-
-# Generate with LoRA style
-image = pipe("A portrait in the trained style").images[0]
-
-# Adjust LoRA strength
-pipe.fuse_lora(lora_scale=0.8)
-
-# Unload LoRA
-pipe.unload_lora_weights()
-```
-
-### 多个 LoRA
-
-```python
-# Load multiple LoRAs
-pipe.load_lora_weights("lora1", adapter_name="style")
-pipe.load_lora_weights("lora2", adapter_name="character")
-
-# Set weights for each
-pipe.set_adapters(["style", "character"], adapter_weights=[0.7, 0.5])
-
-image = pipe("A portrait").images[0]
-```
-
-## 内存优化
-
-### 启用 CPU 卸载
-
-```python
-# Model CPU offload - moves models to CPU when not in use
-pipe.enable_model_cpu_offload()
-
-# Sequential CPU offload - more aggressive, slower
-pipe.enable_sequential_cpu_offload()
-```
-
-### Attention 切片
-
-```python
-# Reduce memory by computing attention in chunks
-pipe.enable_attention_slicing()
-
-# Or specific chunk size
-pipe.enable_attention_slicing("max")
-```
-
-### xFormers 内存高效 Attention
-
-```python
-# Requires xformers package
-pipe.enable_xformers_memory_efficient_attention()
-```
-
-### 大图像的 VAE 切片
-
-```python
-# Decode latents in tiles for large images
-pipe.enable_vae_slicing()
-pipe.enable_vae_tiling()
-```
-
-## 模型变体
-
-### 加载不同精度
-
-```python
-# FP16 (recommended for GPU)
-pipe = DiffusionPipeline.from_pretrained(
-    "model-id",
-    torch_dtype=torch.float16,
-    variant="fp16"
-)
-
-# BF16 (better precision, requires Ampere+ GPU)
-pipe = DiffusionPipeline.from_pretrained(
-    "model-id",
-    torch_dtype=torch.bfloat16
-)
-```
-
-### 加载特定组件
-
-```python
-from diffusers import UNet2DConditionModel, AutoencoderKL
-
-# Load custom VAE
-vae = AutoencoderKL.from_pretrained("stabilityai/sd-vae-ft-mse")
-
-# Use with pipeline
-pipe = DiffusionPipeline.from_pretrained(
-    "stable-diffusion-v1-5/stable-diffusion-v1-5",
-    vae=vae,
-    torch_dtype=torch.float16
-)
-```
-
-## 批量生成
-
-高效生成多张图像：
-
-```python
-# Multiple prompts
-prompts = [
-    "A cat playing piano",
-    "A dog reading a book",
-    "A bird painting a picture"
-]
-
-images = pipe(prompts, num_inference_steps=30).images
-
-# Multiple images per prompt
-images = pipe(
-    "A beautiful sunset",
-    num_images_per_prompt=4,
-    num_inference_steps=30
-).images
-```
-
-## 常见工作流
-
-### 工作流 1：高质量生成
-
-```python
-from diffusers import StableDiffusionXLPipeline, DPMSolverMultistepScheduler
-import torch
-
-# 1. Load SDXL with optimizations
-pipe = StableDiffusionXLPipeline.from_pretrained(
-    "stabilityai/stable-diffusion-xl-base-1.0",
-    torch_dtype=torch.float16,
-    variant="fp16"
-)
-pipe.to("cuda")
-pipe.scheduler = DPMSolverMultistepScheduler.from_config(pipe.scheduler.config)
-pipe.enable_model_cpu_offload()
-
-# 2. Generate with quality settings
-image = pipe(
-    prompt="A majestic lion in the savanna, golden hour lighting, 8k, detailed fur",
-    negative_prompt="blurry, low quality, cartoon, anime, sketch",
-    num_inference_steps=30,
-    guidance_scale=7.5,
-    height=1024,
-    width=1024
-).images[0]
-```
-
-### 工作流 2：快速原型验证
-
-```python
-from diffusers import AutoPipelineForText2Image, LCMScheduler
-import torch
-
-# Use LCM for 4-8 step generation
-pipe = AutoPipelineForText2Image.from_pretrained(
-    "stabilityai/stable-diffusion-xl-base-1.0",
-    torch_dtype=torch.float16
-).to("cuda")
-
-# Load LCM LoRA for fast generation
-pipe.load_lora_weights("latent-consistency/lcm-lora-sdxl")
-pipe.scheduler = LCMScheduler.from_config(pipe.scheduler.config)
-pipe.fuse_lora()
-
-# Generate in ~1 second
-image = pipe(
-    "A beautiful landscape",
-    num_inference_steps=4,
-    guidance_scale=1.0
-).images[0]
-```
-
-## 常见问题
-
-**CUDA 内存不足：**
-```python
-# Enable memory optimizations
-pipe.enable_model_cpu_offload()
-pipe.enable_attention_slicing()
-pipe.enable_vae_slicing()
-
-# Or use lower precision
-pipe = DiffusionPipeline.from_pretrained(model_id, torch_dtype=torch.float16)
-```
-
-**黑色/噪声图像：**
-```python
-# Check VAE configuration
-# Use safety checker bypass if needed
-pipe.safety_checker = None
-
-# Ensure proper dtype consistency
-pipe = pipe.to(dtype=torch.float16)
-```
-
-**生成速度慢：**
-```python
-# Use faster scheduler
-from diffusers import DPMSolverMultistepScheduler
-pipe.scheduler = DPMSolverMultistepScheduler.from_config(pipe.scheduler.config)
-
-# Reduce steps
-image = pipe(prompt, num_inference_steps=20).images[0]
-```
-
-## 参考资料
-
-- **[高级用法](https://github.com/NousResearch/hermes-agent/blob/main/optional-skills/mlops/stable-diffusion/references/advanced-usage.md)** - 自定义 pipeline、微调、部署
-- **[故障排查](https://github.com/NousResearch/hermes-agent/blob/main/optional-skills/mlops/stable-diffusion/references/troubleshooting.md)** - 常见问题与解决方案
-
-## 资源
-
-- **文档**：https://huggingface.co/docs/diffusers
-- **代码仓库**：https://github.com/huggingface/diffusers
-- **模型中心**：https://huggingface.co/models?library=diffusers
-- **Discord**：https://discord.gg/diffusers